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BITSAT 2023 Paper Memory Based

This document contains solutions to 8 math memory-based questions from a BITSAT question paper. The questions cover topics like differential equations, distance between a point and line, definite integrals, and probability. Complete solutions involving relevant formulas, concepts and steps are provided for each question.

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0% found this document useful (0 votes)
58 views19 pages

BITSAT 2023 Paper Memory Based

This document contains solutions to 8 math memory-based questions from a BITSAT question paper. The questions cover topics like differential equations, distance between a point and line, definite integrals, and probability. Complete solutions involving relevant formulas, concepts and steps are provided for each question.

Uploaded by

krishbhatia1503
Copyright
© © All Rights Reserved
We take content rights seriously. If you suspect this is your content, claim it here.
Available Formats
Download as PDF, TXT or read online on Scribd
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BITSAT 2023 Question Paper with Answers PDF

May 22 Shift 1 (Memory-based)

Maths Memory Based Questions

Question 1. General solution of differential equation

Ans. The general solution of a differential equation refers to the set of all possible
solutions that satisfy the equation. The specific form of the general solution
depends on the type and order of the differential equation.

Different types of differential equations have different methods for finding their
general solutions. Here are a few common types of differential equations and
their general solutions:

Ordinary Differential Equations (ODEs):

First-Order Linear ODE: y = C*e^(-∫P(x)dx) + ∫[Q(x)e^(-∫P(x)dx)]dx


Second-Order Homogeneous ODE: y = C1e^(r1x) + C2e^(r2x)
Second-Order Non-Homogeneous ODE (with constant coefficients): y = yh + yp,
where yh is the general solution of the associated homogeneous equation and yp
is a particular solution.
Partial Differential Equations (PDEs):

Separation of Variables: The general solution is found by separating variables


and solving each resulting ordinary differential equation independently.
Fourier Series: The general solution is represented as an infinite series involving
trigonometric functions (sine and cosine) or other orthogonal functions.
It's important to note that the general solution may contain arbitrary constants
(C1, C2, etc.) that are determined by applying initial or boundary conditions to
obtain a particular solution that satisfies specific constraints.
The specific method for finding the general solution of a differential equation
depends on the equation's type, order, and particular characteristics. More
advanced techniques, such as Laplace transforms, Fourier transforms, or
numerical methods, may be required for more complex differential equations.

To obtain the general solution for a specific differential equation, it is necessary to


determine the equation's type and apply the appropriate methods and techniques
accordingly.

Question 2. Dist. of point from line

Ans. To find the distance between a point and a line in a two-dimensional space,
you can use the formula for the distance between a point and a line.

Consider a line defined by the equation Ax + By + C = 0, where A, B, and C are


constants representing the coefficients of the line equation. Let (x0, y0) be the
coordinates of the point.

The formula to calculate the distance between the point and the line is given by:

d = |Ax0 + By0 + C| / √(A^2 + B^2)

where |Ax0 + By0 + C| denotes the absolute value of the expression Ax0 + By0 +
C.

By substituting the appropriate values for A, B, C, x0, and y0, you can calculate
the distance (d) between the point and the line.

It's worth noting that this formula assumes a straight line and Euclidean distance
in a two-dimensional Cartesian coordinate system. If you are working in a
different coordinate system or considering a non-linear curve, the calculation
method may vary.

Question 3. Lpp area of solution region


Question 4. Fog function

Ans. The fog function, also known as a composite function, is a mathematical


operation that combines two functions together. It is denoted by "f ∘ g" and is
read as "f composed with g."

Given two functions, f(x) and g(x), the fog function is defined as:

(f ∘ g)(x) = f(g(x))

In other words, to evaluate the fog function at a specific value of x, you first apply
the function g(x) to x, and then apply the resulting value to the function f(x).

The order of the composition is important. The function g(x) is applied first,
followed by the function f(x). The output of g(x) serves as the input to f(x).

The fog function allows you to combine functions to create more complex
functions and perform calculations in a step-by-step manner. By applying one
function to the output of another, you can create a chain of operations.

Question 5. Definite integration with Wall's formula

Ans. I apologize, but I am not familiar with a specific mathematical formula called
"Wall's formula" related to definite integration. It's possible that you may be
referring to a formula or method that I am not aware of or that is not widely
recognized by that name.

Definite integration typically involves finding the area under a curve between two
given limits. It is often denoted by the integral symbol (∫) and has the form:

∫[a, b] f(x) dx

where "a" and "b" represent the lower and upper limits of integration, "f(x)"
represents the integrand (the function being integrated), and "dx" represents the
differential element.
If you can provide more specific details or clarify the context or formula you are
referring to as "Wall's formula," I would be happy to assist you further.

Question 5. Conditional probability with given P(A) and P(AUB) and asked
to find P(B)
two questions of this type

Ans. If we are given the probability of event A, P(A), and the probability of the
union of events A and B, P(AUB), and we are asked to find the probability of
event B, P(B), we can use the formula for conditional probability.

The conditional probability of event B given event A is denoted as P(B|A) and can
be calculated using the following formula:

P(B|A) = P(A∩B) / P(A)

Where P(A∩B) represents the probability of the intersection of events A and B.

Now, let's rearrange the formula to solve for P(B):

P(B|A) * P(A) = P(A∩B)

P(B) = P(A∩B) / P(A|B)

Given that P(A) and P(AUB) are provided, we can use these probabilities to
calculate P(B) as follows:

P(A∩B) = P(AUB) - P(AUB') (Using the formula for the probability of the
intersection of two events)

P(A|B) = P(A∩B) / P(B) (Using the formula for conditional probability)

Simplifying, we can express P(B) in terms of the given probabilities:

P(B) = P(A∩B) / P(A|B) = (P(AUB) - P(AUB')) / P(A|B)


Therefore, to find P(B), we need to know the conditional probability P(A|B) or
additional information about the relationship between events A and B.

Question 6. (b-a) /2 thing of Definite integrals ( integrate from a to b)

Ans. The expression (b - a) / 2 represents half the difference between two


values, b and a. In the context of definite integrals, the expression (b - a) / 2 often
arises when using certain numerical integration methods, such as the midpoint
rule.

The midpoint rule is a simple numerical approximation method for evaluating


definite integrals. It divides the interval [a, b] into subintervals and approximates
the area under the curve by treating each subinterval as a rectangle with a height
determined by the function's value at the midpoint of the subinterval.

To apply the midpoint rule, you first divide the interval [a, b] into n equal
subintervals. The width of each subinterval is given by Δx = (b - a) / n. Then, you
evaluate the function at the midpoint of each subinterval and multiply it by the
width of the subinterval. Finally, you sum up all these approximations to get an
estimate of the definite integral.

Mathematically, the approximation of the definite integral using the midpoint rule
can be written as:

∫[a, b] f(x) dx ≈ Δx * (f(a + Δx/2) + f(a + 3Δx/2) + f(a + 5Δx/2) + ... + f(b - Δx/2))

where Δx = (b - a) / n and n is the number of subintervals.

So, (b - a) / 2 appears as half the width of each subinterval in the midpoint rule,
as it calculates the midpoint for each subinterval.

Question 7. INTEGRATION SIN4X.COS2X INTEGRATION LIMIT PI/2

Ans. To evaluate the definite integral of sin^4(x) * cos(2x) with integration limits
from π/2, we can use the power-reducing identities for sine and double angle
identity for cosine. Let's break down the steps:
Apply the power-reducing identity for sine:
sin^4(x) = (1/2)(1 - cos(2x))^2

Apply the double angle identity for cosine:


cos(2x) = 2cos^2(x) - 1

Now we have:
∫(π/2) sin^4(x) * cos(2x) dx
= ∫(π/2) (1/2)(1 - cos(2x))^2 * (2cos^2(x) - 1) dx

Expanding and simplifying:


= (1/2) * ∫(π/2) (1 - 2cos(2x) + cos^2(2x)) * (2cos^2(x) - 1) dx

Distribute and simplify further:


= (1/2) * ∫(π/2) (2cos^2(x) - 4cos^3(2x) + 2cos^4(2x) - 2cos^2(x) + cos^2(2x)) dx

Combining like terms:


= (1/2) * ∫(π/2) (cos^2(x) - 4cos^3(2x) + 2cos^4(2x) + cos^2(2x)) dx

Now we have separate terms that we can integrate individually:


Term 1: ∫(π/2) cos^2(x) dx
Term 2: -4 * ∫(π/2) cos^3(2x) dx
Term 3: 2 * ∫(π/2) cos^4(2x) dx
Term 4: ∫(π/2) cos^2(2x) dx

Each of these terms can be integrated using basic integration techniques. Once
integrated, you can substitute the upper limit π/2 and subtract the result with the
lower limit to obtain the definite integral.

Question 8. Binomial standard formula CO + C1 +…....nCn

Ans. The binomial standard formula, also known as the binomial theorem,
expresses the expansion of a binomial raised to a positive integer power. It can
be stated as follows:

(x + y)^n = C(n, 0) * x^n * y^0 + C(n, 1) * x^(n-1) * y^1 + C(n, 2) * x^(n-2) * y^2 +
... + C(n, n) * x^0 * y^n
In this formula, (x + y) represents a binomial expression with x and y as
variables, n is a positive integer exponent, and C(n, k) represents the binomial
coefficient. The binomial coefficient C(n, k) represents the number of ways to
choose k items from a set of n items and is given by the formula:

C(n, k) = n! / (k! * (n-k)!)

Now, let's rewrite the formula using the notation you provided:

CO + C1 + ... + nCn

This notation represents the sum of binomial coefficients from C0 to Cn. It can be
written using the binomial standard formula as follows:

C0 + C1 + ... + Cn = (x + y)^n

In other words, the sum of all binomial coefficients from C0 to Cn is equal to the
binomial expansion of (x + y)^n.

Question 9. how many cases are possible if a batsman can hit 0, 1,2,3,4,6
on 6 balls to make total score 30

Ans. To determine the number of cases in which a batsman can score a total of
30 runs by hitting 0, 1, 2, 3, 4, or 6 on each of the 6 balls, we can use
combinatorics. We need to find the number of combinations of scores that sum
up to 30.

Let's break it down:

We have 6 balls to consider, and each ball can have 6 possible scores (0, 1, 2, 3,
4, or 6).

We need to find the number of combinations of these scores that sum up to 30.

We can solve this problem using a technique called "stars and bars" or "balls and
urns."
Consider the 6 balls as 6 "stars." We need to distribute these stars among the
scores (0, 1, 2, 3, 4, and 6) to get a total of 30 runs. To represent the distribution,
we can use "bars" or "urns" between the scores.

For example, if we have the distribution: 6 | 4 | 6 | 6 | 4 | 4, it means the batsman


hit a 6 on the first ball, a 4 on the second ball, a 6 on the third ball, a 6 on the
fourth ball, a 4 on the fifth ball, and a 4 on the sixth ball.

To find the number of cases, we need to find the number of ways to arrange the 6
stars and 5 bars. The bars divide the stars into different scores. We can calculate
this using the concept of combinations.

The number of cases is given by:

Number of cases = C(stars + bars, bars)

In this case, we have 6 stars and 5 bars, so:

Number of cases = C(6 + 5, 5) = C(11, 5)

Using the binomial coefficient formula:

C(n, k) = n! / (k! * (n-k)!)

we can calculate the number of cases as:

C(11, 5) = 11! / (5! * (11-5)!) = 11! / (5! * 6!)

Calculating this expression:

Number of cases = 462

Therefore, there are 462 possible cases in which a batsman can score a total of
30 runs by hitting 0, 1, 2, 3, 4, or 6 on each of the 6 balls.

Question 10. 3j +5z,1-4]+z, 2i+ 5j-kz they are collinear find k


Ans. To determine if the vectors 3j + 5z, [1-4]+z, and 2i + 5j - kz are collinear, we
need to check if they are scalar multiples of each other. This means that one
vector can be obtained by multiplying another vector by a scalar constant.

Let's compare the components of the given vectors:

Vector 1: 3j + 5z
Vector 2: [1-4]+z
Vector 3: 2i + 5j - kz

Comparing the x-components, we have 0 in Vector 1 and 1 in Vector 3. They are


not equal, so the vectors cannot be scalar multiples of each other.

Therefore, the vectors 3j + 5z, [1-4]+z, and 2i + 5j - kz are not collinear.

Since the vectors are not collinear, we do not need to find the value of k.

Chemistry Memory Based Questions

Question 1. pH order of 0. 1M OF sodium acetate, nh4cl and nacl

Ans. To determine the pH order of three solutions: 0.1 M sodium acetate, NH4Cl,
and NaCl, we need to consider the nature of the dissolved compounds and their
effect on the pH of the solution.

Sodium Acetate (CH3COONa):


Sodium acetate is a salt of a weak acid (acetic acid) and a strong base (sodium
hydroxide). When sodium acetate dissolves in water, it undergoes hydrolysis,
which results in the formation of acetate ions (CH3COO-) and a small amount of
hydroxide ions (OH-). Acetate ions being the conjugate base of a weak acid, are
relatively basic.
As a result, the presence of sodium acetate in water will slightly increase the pH
of the solution. Therefore, the pH of the 0.1 M sodium acetate solution will be
slightly basic, but the extent of the increase will depend on the concentration.
NH4Cl:
NH4Cl is a salt of a weak base (ammonia) and a strong acid (hydrochloric acid).
When NH4Cl dissolves in water, it dissociates into ammonium ions (NH4+) and
chloride ions (Cl-). Ammonium ions are the conjugate acid of a weak base and
are relatively acidic.
The presence of NH4Cl in water will slightly decrease the pH of the solution.
Therefore, the pH of the NH4Cl solution will be slightly acidic, but again, the
extent of the decrease will depend on the concentration.

NaCl:
NaCl is a salt of a strong acid (hydrochloric acid) and a strong base (sodium
hydroxide). When NaCl dissolves in water, it dissociates into sodium ions (Na+)
and chloride ions (Cl-). Neither sodium nor chloride ions significantly affect the
pH of the solution.
Therefore, the pH of a NaCl solution will be neutral, around 7.

In summary, the pH order of the given solutions from most acidic to most basic
would be:
NH4Cl < NaCl < 0.1 M sodium acetate.
Question 2. which has least basic conjugate base options hcl,hbr,oxalic
acid and acetic acid

Ans. To determine which compound has the least basic conjugate base among
HCl, HBr, oxalic acid, and acetic acid, we need to consider the acidity of their
respective conjugate bases.

The conjugate base of an acid is formed when the acid donates a proton (H+).
The stronger the acid, the weaker its conjugate base.

Let's analyze each compound:

HCl (hydrochloric acid):


The conjugate base of HCl is chloride ion (Cl-). Chloride ion is very weakly basic.

HBr (hydrobromic acid):


The conjugate base of HBr is bromide ion (Br-). Bromide ion is also very weakly
basic.

Oxalic acid:
Oxalic acid has two acidic protons, and its conjugate base is oxalate ion
(C2O4^2-). Oxalate ion is a weak base but stronger than chloride or bromide
ions.

Acetic acid:
Acetic acid has a weak acidity, and its conjugate base is acetate ion (CH3COO-).
Acetate ion is a relatively weak base but stronger than chloride, bromide, or
oxalate ions.

Therefore, among the given compounds, oxalic acid has the least basic
conjugate base.

In summary, the order of least basic conjugate bases is:


HCl < HBr < acetic acid < oxalic acid

Question 3. molecular orbital config of N2

Ans. The molecular orbital configuration of N2 (nitrogen gas) can be determined


by combining the atomic orbitals of two nitrogen atoms. Nitrogen has atomic
number 7, so each nitrogen atom has 7 electrons.

The molecular orbital diagram for N2 can be represented as follows:

σ2s (bonding)
σ2s* (anti-bonding)
σ2p (bonding)
σ2p* (anti-bonding)
π2p (bonding)
π2p* (anti-bonding)
In this diagram, the σ2s and σ2p orbitals are lower in energy and bonding
orbitals, while the σ2s* and σ2p* are higher in energy and anti-bonding orbitals.
The π2p and π2p* orbitals are degenerate and form bonding and anti-bonding
orbitals.

For N2, we need to fill these molecular orbitals with a total of 14 electrons (7
electrons from each nitrogen atom). We start by filling the orbitals from lowest to
highest energy, following the Pauli exclusion principle and Hund's rule.

The electron configuration for N2 is as follows:

σ2s: ↑↓
σ2s*: (empty)
σ2p: ↑↓
σ2p*: (empty)
π2p: ↑↑↓↓
π2p*: (empty)

The molecular orbital configuration shows that N2 has a triple bond. The π2p
bonding orbitals are fully occupied with 4 electrons, resulting in a strong covalent
bond between the nitrogen atoms. This triple bond is responsible for the stability
and strength of the nitrogen molecule.

It's important to note that the diagram and configuration provided here are
simplified representations of the molecular orbitals in N2. The actual molecular
orbital energies and electron distributions can be more complex, but this serves
as a basic understanding of the molecular orbital configuration of N2

Question 4. is NaOH a standard titrant

Ans. Yes, NaOH (sodium hydroxide) is commonly used as a standard titrant in


acid-base titrations. In acid-base titrations, a solution of known concentration (the
titrant) is used to determine the concentration of an unknown solution (the
analyte) by measuring the volume of titrant required to reach a specific endpoint.

NaOH is considered a standard titrant because it can be prepared with high


accuracy and precision, and its concentration can be determined through primary
standardization methods. A primary standard is a highly pure compound that can
be accurately weighed and dissolved to prepare a solution of known
concentration.

In the case of NaOH, its concentration can be determined by titration against a


standardized acid solution (such as hydrochloric acid) using a suitable indicator
or pH meter to determine the endpoint of the reaction.

Once the concentration of the NaOH solution is accurately determined, it can be


used as a standard titrant in various acid-base titrations to determine the
concentration of acidic analytes.

Question 5. miller indices(solid state)

Ans. Miller indices are a notation system used in solid-state physics and
crystallography to describe the orientation and planes of crystal lattice structures.
They are denoted by three integers (hkl), known as Miller indices, which
represent the reciprocals of the fractional intercepts made by a plane on the
crystallographic axes.

The steps to determine the Miller indices of a plane in a crystal lattice are as
follows:

Identify the intercepts: Determine the points at which the plane intersects the
crystallographic axes (x, y, z).

Take the reciprocals: Take the reciprocals of the fractional intercepts to obtain
integers. If there is no intercept along a particular axis, assign a value of infinity
(∞) or zero (0).

Simplify the indices: If the indices have any common factors, divide all of them by
the greatest common factor to simplify the Miller indices.

Enclose the indices in parentheses: Write the simplified indices as (hkl), where h,
k, and l represent the integers obtained in steps 2 and 3.
For example, if a plane intersects the x-axis at 1/2, the y-axis at 1/3, and does
not intersect the z-axis, the Miller indices would be:

h = 2 (reciprocal of 1/2)
k = 3 (reciprocal of 1/3)
l = 0 (no intercept along the z-axis)

After simplifying, the Miller indices would be written as (2 3 0).

Miller indices are used to describe crystallographic planes, directions, and other
lattice features in solid-state materials. They provide a compact and standardized
way of representing the orientation and symmetry of crystal structures.

Question 6. when calcium is titrated with EDTA what is the ph of buffer

Ans. The pH of the buffer used in the titration of calcium with EDTA
(ethylenediaminetetraacetic acid) depends on the specific composition of the
buffer solution.

In general, when titrating calcium with EDTA, a common choice for the buffer is
ammonia-ammonium chloride buffer, commonly known as "NH4Cl/NH3 buffer" or
"ammonium buffer." This buffer system consists of a mixture of ammonium
chloride (NH4Cl) and ammonium hydroxide (NH3) in water.

The pH of the ammonium buffer can be adjusted by controlling the relative


concentrations of NH4Cl and NH3. The buffer capacity is typically achieved
around pH 9-10, which is suitable for complexation reactions involving calcium
and EDTA.

The actual pH of the buffer solution used in the titration would depend on the
specific concentrations of NH4Cl and NH3, as well as the temperature and ionic
strength of the solution. Therefore, to determine the precise pH value of the
buffer, the specific composition and concentrations of the buffer components
need to be known or specified.

It's important to note that maintaining a consistent and appropriate pH during the
titration is crucial for accurate results, as the complex formation between calcium
and EDTA is pH-dependent. Therefore, careful consideration and adjustment of
the buffer pH are necessary to ensure the success of the calcium titration with
EDTA.

Physics Memory Based Questions

Question 1. If a ring and disc are allowed to come drown through inclined
plane, then which will come down first

Ans. The ring and the disc will have different rates of descent when allowed to
roll down an inclined plane due to their different moments of inertia and
distribution of mass.

When a ring rolls down an inclined plane, its moment of inertia is higher
compared to a disc of the same mass and radius. This is because the mass of
the ring is distributed farther from its rotational axis, resulting in a larger moment
of inertia. Due to its larger moment of inertia, the ring tends to resist rotational
motion more than the disc.

As a result, the disc will come down the inclined plane first, followed by the ring.
The disc's lower moment of inertia allows it to gain rotational motion and roll
down the incline more easily than the ring.

Therefore, the disc will descend first, while the ring will take longer to roll down
the inclined plane.

Question 2. Ek work nikalna tha f diya tha vector form mai s diya tha ek
point sai duse point tak

Ans. To calculate the work done (W) given a force vector F and displacement
vector s, you can use the dot product between the force vector and the
displacement vector. The formula for calculating work in vector form is:

W=F·s
where · represents the dot product.

The dot product of two vectors is calculated by multiplying the corresponding


components of the vectors and taking their sum. If F = (Fx, Fy, Fz) and s = (sx,
sy, sz), then the dot product F · s is given by:

F · s = Fx * sx + Fy * sy + Fz * sz

So, to calculate the work done, you need the components of the force vector F
and the displacement vector s. Multiply the corresponding components and sum
them up to obtain the scalar value of the work done.

Question 3. Field due to infinite sheet(gave 2 sheets with same charge


density and based net field outside the region common

Ans.
When considering two infinite sheets with the same charge density and opposite
charges, the net electric field outside the region common to both sheets will be
zero.

This can be understood by considering the symmetry of the system. Since the
charge densities on both sheets are the same but have opposite signs, the
electric fields produced by each sheet will cancel each other out in the region
outside where they overlap.

However, it's important to note that the electric field due to each individual sheet
will still exist and contribute to the net field within their respective regions. The
cancellation of the electric fields only occurs in the overlapping region outside
both sheets.

Inside the region common to both sheets, the net electric field will depend on the
distance from each sheet and their respective charge densities. To calculate the
exact electric field within that region, the geometry and specific values of the
charge densities would need to be known.

Question 4. Gave ratio of volume submerged in 2 diff liquids and they


asked ratio of densities of liquid,
Ans.
If the ratio of volumes submerged in two different liquids is given, we can use
Archimedes' principle to determine the ratio of their densities.

According to Archimedes' principle, the buoyant force experienced by an object


submerged in a fluid is equal to the weight of the fluid displaced by the object.
The buoyant force can be expressed as:

F_buoyant = ρ_fluid * V_submerged * g

where ρ_fluid is the density of the fluid, V_submerged is the volume of the object
submerged in the fluid, and g is the acceleration due to gravity.

If we have two different liquids, let's say Liquid 1 and Liquid 2, and their
respective volumes submerged are V_submerged_1 and V_submerged_2, we
can set up the following equation:

ρ_1 * V_submerged_1 * g = ρ_2 * V_submerged_2 * g

Canceling out the common factors, we find:

ρ_1 / ρ_2 = V_submerged_2 / V_submerged_1

Therefore, the ratio of the densities of the two liquids (ρ_1 / ρ_2) is equal to the
ratio of the volumes submerged (V_submerged_2 / V_submerged_1).

By knowing the ratio of volumes submerged, you can directly determine the ratio
of densities of the two liquids.

Question 5. Block is thrown on inclined plane with friction and asked if


velocity with which it will come back be treated or less

Ans. When a block is thrown on an inclined plane with friction, the velocity with
which it will come back (assuming it returns along the same path) will be less
than the velocity at which it was thrown.
Several factors contribute to this decrease in velocity:

Friction: The presence of friction between the block and the inclined plane will
oppose its motion. Friction converts some of the kinetic energy into heat energy,
causing a loss of mechanical energy. Consequently, the block will lose speed as
it moves up the incline.

Work done against gravity: As the block moves up the incline, it works against
gravity. This work done by the block results in a loss of kinetic energy and a
decrease in velocity.

Losses due to air resistance: In real-world scenarios, air resistance can also play
a role in reducing the velocity of the block. Air resistance acts opposite to the
direction of motion and can further diminish the velocity of the block as it moves
against it.

Overall, due to the combined effects of friction, work done against gravity, and air
resistance (if applicable), the velocity with which the block will come back up the
inclined plane will be less than the velocity at which it was thrown.

Question 6. work done by tension in circular motion

Ans. In circular motion, the tension force does work on an object moving in a
circular path. The work done by tension can be determined by considering the
change in kinetic energy of the object.

When an object is moving in a circular path at a constant speed, its velocity is


constantly changing direction. This change in direction requires a force to act
towards the center of the circular path, which is provided by the tension force in
this case.

Since the tension force is perpendicular to the displacement of the object, the
work done by tension is zero. This is because the dot product of the tension force
and displacement vector is zero when the angle between them is 90 degrees.

Mathematically, the work done by tension (W_tension) can be expressed as:


W_tension = T * s * cosθ

where T is the tension force, s is the displacement of the object, and θ is the
angle between the tension force and displacement vector. Since cos(90°) = 0, the
work done by tension becomes zero.

Therefore, in circular motion, the tension force does not do any work on the
object because the displacement is perpendicular to the tension force. The work
done by tension is zero.

Other Memory Based Questions

Question 1. -3,7,31,2,145; pattern completion

Ans. Looking at the given sequence of numbers: -3, 7, 31, 2, 145, it is not
immediately evident what the pattern or rule is for completing the sequence.
However, let's analyze the numbers further to see if any patterns emerge.

If we calculate the differences between consecutive terms, we get:

7 - (-3) = 10
31 - 7 = 24
2 - 31 = -29
145 - 2 = 143

The differences between the terms are not following a consistent pattern. It
appears that there might not be a simple mathematical pattern or rule governing
the sequence. It's also possible that there may be a pattern or rule that is not
immediately obvious or that additional terms are required to discern the pattern
accurately.

Without further information or context, it is difficult to determine the exact pattern


or complete the sequence confidently.

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